Все модели
Полный список отсортирован по популярности на Replicate.
Lucid Sonic Dreams syncs StyleGAN XL -generated visuals to music
Juggernaut Aftermath Model with original TRCVAE (Text2Img, Img2Img and Inpainting)
Anime Aesthetic Predict
🕹️FramePack: video diffusion that feels like image diffusion🎥
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDXL images with an image prompt
AI that transforms sketches into realistic images. Upload your drawing and describe it in the prompt. You can also adjust the ControlNet parameters and scale the image to a higher resolution for better results
Flux LoRA for creating geometrical styled illustrations. Use "illustration in the style of NWGMTRCPOPV01" to trigger the model.
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDv1.5 images with an image prompt
DiT-based 13b video generation model, creating 30fps video
SDXL, but faster
Upscale | Enhancer | Ultra-Resolution | Restoration |
Recently, this variant of Code Llama called WizardCoder 34B has emerged and achieved better scores than GPT4 on Human Eval
Classifies Pokémon Yellow game screens for automated gameplay (Battle, Menu, Overworld, Dialogue)
Trigger phrase: surreal style
Japanese-specific latent text-to-image diffusion model
Fast video2video with a latent consistency model
Training-Free Text-to-Image Generation
DeepSeek LLM, an advanced language model comprising 67 billion parameters. Trained from scratch on a vast dataset of 2 trillion tokens in both English and Chinese
Stable Diffusion fined-tuned on frames from Monkey Island 1 and 2
Training-free Controllable Text-to-Video Generation
Dolphin is uncensored. I have filtered the dataset to remove alignment and bias. This makes the model more compliant.
Just an experiment. Nothing new here.
App Icons Generator V1 (DreamBooth Model)
Generate 5s and 10s videos in 1080p resolution at 30fps
Deep Flexible Structure-preserving Image Smoothing
a powerful and competitive model like Midjourney v6 and DALL-E 3 but Open and Decentralized
Photo to Anime – Stylized conversion that turns photos into crisp cel-shaded anime frames using the Photo-to-Anime LoRA.
Generate a video transitioning from one image to another using SEINE model
Stable Diffusion v2-1-unclip Model
Stable Diffusion 3.5 Large - LoRA Explorer
A LoRA for Flux.1 Dev to re-create really bad and trippy CGI from the 90s.
ControlNet for FLUX.1-dev model jointly released by InstantX and Shakker Labs
Anydoor: zero-shot object-level image customization
A Visual Language Model for GUI Agents
Qwen image edit fast
Higgs Audio v2, a powerful text-to-speech audio foundation model that excels in expressive audio generation
🖼️ Super fast 1.5B Image Captioning/VQA Multimodal LLM (Image-to-Text) 🖋️
Test deployment of OuteTTS 500M
Generate haiku. A tiny model for testing out Replicate.
Prompt-to-prompt image editing with cross-attention control
Generate high-quality images faster with Latent Consistency Models (LCM), a novel approach that distills the original model, reducing the steps required from 25-50 to just 4-8 in Stable Diffusion (SDXL) image generation.
Make a very realistic looking real-world AI video
Extract structured data from receipt images using Donut 🍩 (Document Understanding Transformer)